stable-diffusion-image-generation
Category: multimodal-and-emerging · License: MIT · Version: 1.0.0
Tags: Image Generation, Stable Diffusion, Diffusers, Text-to-Image, Multimodal, Computer Vision
When to use it
Generates images with Stable Diffusion models (SD 1.5, SDXL, SD 3.0, Flux) through the HuggingFace Diffusers library, covering text-to-image, image-to-image, inpainting, outpainting, ControlNet conditioning, LoRA adapters, scheduler swapping, and GPU memory optimization. Use when generating images from text prompts, transforming or restyling an existing image, filling masked regions of an image, adding spatial control from edges, poses, or depth maps, loading LoRA style adapters, or fixing out-of-memory and black-image errors in a Diffusers pipeline. Not for API-only generation without a GPU, where DALL-E 3 fits better.Example requests
Ask your agent in plain language, for example:- “How do I make SDXL fit on my 8GB card? I keep getting CUDA out of memory in diffusers.”
- “I’m running StableDiffusionXLPipeline.from_pretrained(‘stabilityai/stable-diffusion-xl-base-1.0’, torch_dtype=torch.float16) on an RTX 3090 and every image comes out as pure black. A warning says ‘invalid value encountered in cast’. How do I fix it? I’m making figures for a methods paper on histology slide augmentation.”
- “I have a set of grayscale microscope photos of cell cultures. I want to turn each one into a colorized, painted look while keeping the cell layout and shapes the same. I have a local GPU. What’s a good way to do this?”
Install
/plugin marketplace add KalarisLabs/research-agent-skills and install the whole multimodal-and-emerging category with /plugin install multimodal-and-emerging@research-agent-skills.